Challenge: Existing methods for inference are expensive and lack spatial redundancy . Discrete Diffusion Language Models are a promising paradigm for multimodal generation .
Approach: They propose a locality-aware dynamic rescue method that exploits spatial Markov property of images.
Outcome: The proposed method achieves an approximate 4 speedup over baselines on four text-to-image generation benchmarks.

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LaDiC: Are Diffusion Models Really Inferior to Autoregressive Counterparts for Image-to-Text Generation? (2024.naacl-long)

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Challenge: Existing models for text-to-image generation have been underperforming in image-totext generation tasks.
Approach: They propose a framework that uses a split BERT to create a dedicated latent space for captions and integrates a regularization module to manage varying text lengths.
Outcome: The proposed framework achieves state-of-the-art performance on the MS COCO dataset with 38.2 BLEU@4 and 126.2 CIDEr .
Rapid Diffusion: Building Domain-Specific Text-to-Image Synthesizers with Fast Inference Speed (2023.acl-industry)

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Challenge: Text-to-Image Synthesis (TIS) aims to generate images based on textual inputs . but, current diffusion-based models lack entity knowledge and low inference speed .
Approach: They propose a framework for training and deploying latent diffusion models with rich entity knowledge injected and optimized networks.
Outcome: The proposed framework improves image quality and inference speed and can be used in industrial applications.
LAD: LoRA-Adapted Diffusion (2025.emnlp-demos)

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Challenge: Autoregressive models dominate text generation but suffer from left-to-right decoding constraints that limit efficiency and bidirectional reasoning.
Approach: They propose a framework for non-autoregressive generation that adapts LLaMA models for iterative, bidirectional sequence refinement using LoRA adapters.
Outcome: The proposed framework adapts LLaMA models for iterative, bidirectional sequence refinement using LoRA adapters.
Can Diffusion Model Achieve Better Performance in Text Generation ? Bridging the Gap between Training and Inference ! (2023.findings-acl)

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Challenge: Existing models for text generation use a discrete data embedding module to map the data into the continuous space.
Approach: They propose two methods to bridge the gap between training and inference by mapping the discrete text into the continuous space.
Outcome: The proposed methods can achieve 100 200 speedup with better performance on 6 generation tasks.
DiffuSeq-v2: Bridging Discrete and Continuous Text Spaces for Accelerated Seq2Seq Diffusion Models (2023.findings-emnlp)

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Challenge: Existing approaches to text generation use discrete text within a continuous diffusion space, which incurs substantial computational overhead during training and results in slower sampling speeds.
Approach: They propose a soft absorbing state that facilitates diffusion models in learning to reconstruct discrete mutations based on the underlying Gaussian space.
Outcome: The proposed method accelerates training convergence by 4x and generates samples of similar quality 800x faster, rendering it closer to practical application.
From AR to Diffusion: Efficiently Adapting Large Language Models with Strictly Causal and Elastic Horizons (2026.acl-long)

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Challenge: Autoregressive (AR) models rely on bidirectional attention, creating a structural mismatch with pre-trained Autoregression models.
Approach: They propose a framework that efficiently adapts autoregressive (AR) models to the diffusion paradigm.
Outcome: The proposed framework reduces training costs by orders of magnitude while maintaining state-of-the-art performance.
Few-shot Temporal Pruning Accelerates Diffusion Models for Text Generation (2024.lrec-main)

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Challenge: Existing acceleration methods for text generation ignore the importance of the distribution of sampling steps, resulting in slow sampling rates.
Approach: They propose a technique to accelerate diffusion models for text generation without additional training by using a Bayesian optimization approach.
Outcome: The proposed technique achieves 400x acceleration even with minimal sampling steps after down to less than 1 minute of optimization yielding a competitive performance even with minimum sampling steps.
Beyond Fully Random Masking: Attention-Guided Denoising and Optimization for Diffusion Language Models (2026.acl-long)

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Challenge: Existing methods for full-attention dLLMs rely on random masking strategies that overlook intrinsic token dependencies.
Approach: They propose an attention-guided denoising and optimization framework that aligns training and optimization with attention-derived dependencies.
Outcome: The proposed framework outperforms state-of-the-art methods on mathematical and coding benchmarks.
LoPT: Lossless Parallel Tokenization Acceleration for Long Context Inference of Large Language Model (2026.acl-long)

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Challenge: Existing parallel tokenization methods suffer from inconsistent results due to boundary artifacts that occur after merging.
Approach: They propose a Lossless Parallel Tokenization framework that ensures output identical to standard sequential tokenization.
Outcome: The proposed method achieves significant speedup while guaranteeing lossless tokenization.
DOS: Dependency-Oriented Sampler for Masked Diffusion Language Models (2026.findings-acl)

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Challenge: Existing decoding strategies for pre-trained MDLMs rely on token-level uncertainty criteria, while largely overlooking sequence-level information and inter-token dependencies.
Approach: They propose a training-free decoding strategy that leverages inter-token dependencies to inform token updates during generation.
Outcome: Empirical results show that the proposed approach consistently achieves superior performance on both code generation and mathematical reasoning tasks.

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